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Determining damping terms in fractional wave equations

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Published 1 June 2022 © 2022 The Author(s). Published by IOP Publishing Ltd
, , Citation Barbara Kaltenbacher and William Rundell 2022 Inverse Problems 38 075004 DOI 10.1088/1361-6420/ac6b31

0266-5611/38/7/075004

Abstract

This paper deals with the inverse problem of recovering an arbitrary number of fractional damping terms in a wave equation. We develop several approaches on uniqueness and reconstruction, some of them relying on Tauberian theorems that provide relations between the asymptotic behaviour of solutions in time and Laplace domains. The possibility of additionally recovering space-dependent coefficients or initial data is discussed. The resulting methods for reconstructing coefficients and fractional orders in these terms are tested numerically. In addition, we provide an analysis of the forward problem consisting of a multiterm fractional wave equation.

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1. Introduction

Wave phenomena like seismic waves or ultrasound propagation are known to exhibit power law frequency-dependent damping, which can be modelled by time fractional derivative terms in the corresponding pdes. As a matter of fact, the adequate representation of attenuation with power law frequency dependence was perhaps the first application of fractional time derivative concepts and among the driving forces for their development, see Caputo [4] and Caputo and Mainardi [5].Fractionally damped wave equations arise by combining classical balance laws with fractional constitutive relations. More precisely, consider the equation of motion (resulting from a balance of forces) and the representation of the mechanical strain as the symmetric gradient of the mechanical displacement

Equation (1)

where ρ is the mass density, $\vec{u}(x,t)$ the vector of displacements and σ(x, t) the stress tensor. The crucial constitutive equation is now the material law that relates stress σ and strain ɛ. Some often encountered instances are the fractional Newton (also called Scott–Blair model), Voigt, Maxwell, and Zener models

For a survey on fractionally damped acoustic wave equations we refer to [3]. Driven by applications that are not covered by these models, a general model class has been defined by

Equation (2)

with the normalisation d0 = 1, where 0 ⩽ γ0 < γ1 ⋯ ⩽ γJ and 0 ⩽ α0 < α1 ⋯ ⩽ αK cf, e.g., [1, section 3.1.1], [2, 33]. Combination of (2) with (1) leads to a general wave equation, which for simplicity of exposition in the scalar case of one-dimensional mechanics or in acoustics reads as

In this more complicated situation, direct reconstruction of the damping terms from experimental data is not possible any more. It is the aim of this paper to investigate reconstruction of the orders αk , γj and the coefficients bk , dj from time trace measurements of u, which plays the role of a displacement in mechanics or of a pressure in acoustics. In the latter application case, on which we mainly focus here, certain imaging modalities additionally require reconstruction of an unknown sound speed c(x) in b0 = c2 (ultrasound tomography), and/or an initial condition u0(x) (equivalently, of the spatially varying part f of the source r(x, t) = σ(t)f(x); photoacoustic tomography). Besides investigating these tasks, inspired by recent work by Jin and Kian [18] in the context of diffusion models, we will also study the question of identifiability of the damping terms in an unknown medium, that is, e.g., with unknown c(x), without aiming at its reconstruction. While to the best of our knowledge, this is the first work on recovering multiple fractional damping in wave type equations, much more literature is available for anomalous diffusion pdes. Besides the reference [18], we also refer to, e.g., [17, 24, 25, 34, 36].

We will present results on uniqueness and on reconstruction algorithms in each of three paradigms depending on the availability of measured data. These are: a full time trace u(x0, t), t > 0 for some fixed point x0; extreme large time measurements; and extremely small time measurements. The main tool used in the latter two cases is a Tauberian theorem that relates the asymptotic behaviour of the solution in time to the structure of the representation of the solution in the Laplace transformed formulation that contains specific algebraic information on the fractional operator.

The remainder of this paper is organised as follows. In subsections 1.1 and 1.2 we introduce the models under consideration and state the inverse problem. Section 2 is devoted to the forward problem: we provide an analysis of the general midterm fractional wave equation and describe the numerical forward solver used in our reconstructions. In section 3 we prove uniqueness results based on different paradigms: using the Weyl estimate on decay of eigenvalues together with smoothness of the excitation, applying single mode excitations and exploiting information on poles and residues in the Laplace domain. Section 4 develops reconstruction methods in the three scenarios indicated above: full time (that is, Laplace domain), large time and small time observations. Correspondingly, in section 5 we provide numerical reconstructions along with a discussion of the approaches.

1.1. Models

We consider the second order in time damped wave equation

Equation (3)

where $\mathcal{A}=-\mathbb{L}$ on ${\Omega}\subseteq {\mathbb{R}}^{d}$, equipped with homogeneous Dirichlet/Neumann/impedance boundary conditions, for an elliptic differential operator $\mathbb{L}$ and ${\alpha }_{k}\in \left(0,1\right]$, ${\beta }_{k}\in \left(\frac{1}{2},1\right]$.

Throughout this paper, we denote by ${\partial }_{t}^{\alpha }$ the (partial) Caputo–Djrbashian fractional time derivative of order α ∈ (n − 1, n) with $n\in \mathbb{N}$ by

where ${\partial }_{t}^{n}$ denotes the nth integer order partial time derivative and for γ ∈ (0, 1), and ${I}_{t}^{\gamma }$ is the Abel integral operator defined by ${I}_{t}^{\gamma }[v](t)=\frac{1}{{\Gamma}(\gamma )}{\int }_{0}^{t}\frac{v(s)}{{(t-s)}^{1-\gamma }}\,\mathrm{d}s.$ For details on fractional differentiation and subdiffusion equations, we refer to, e.g., [8, 9, 27, 31, 32]. See also the tutorial paper on inverse problems for anomalous diffusion processes [19].

We assume that the differentiation orders with respect to time are distinct, that is

Equation (4)

a property that is crucial for distinguishing the different asymptotic terms in the solution u and its Laplace transform.

More generally than (3), in most of this paper we also take higher than second order time derivatives into account

Equation (5)

Note that the usual models from [21] are contained in (5), which is actually an extension of the general model from, e.g., the book [1] in the sense that it also allows for fractional powers of the operator $\mathcal{A}$.

Typically, in acoustics we just have $\mathbb{L}={\triangle}$, and then the operator $\mathcal{A}$ is known. To take into account a (possibly unknown) spatially varying speed of sound c(x), one can instead consider

Equation (6)

with ${\mathcal{A}}_{c}=-c{(x)}^{2}{\triangle}$; in order to get a selfadjoint operator ${\mathcal{A}}_{c}$ we then use the weighted L2 inner product with weight function $\frac{1}{{c}^{2}}$.

Note that in order to assume (4) without loss of generality, so allowing for all possible combinations of αk , βk , we would have to consider

Equation (7)

or, including higher than second time derivatives such as in (5)

Equation (8)

However, (7) and (8) might be viewed as over-parameterised models and they actually lead to difficulties in proving uniqueness of solutions.

Important special cases of (6) are the Caputo–Wismer–Kelvin–Chen–Holm

Equation (9)

and the fractional Zener fz model

Equation (10)

These pde models will be considered on a time interval t ∈ (0, T) and driven by initial conditions and/or a separable source term

Equation (11)

Note that boundary conditions are already incorporated into the operator $\mathcal{A}$ or ${\mathcal{A}}_{c}$, respectively.

1.2. Inverse problems

As in [18], we assume that data has been obtained from a quite general experimental setup in which the pde (5) with its unknown coefficients remains the same, while both excitation (11) and observation location may vary between the experiments. We number the experiments with an index i to denote the corresponding excitation by (u0,i , u1,i , fi ), and the corresponding observation operator by Bi . Given I experiments, we thus assume to have the following I time trace observations for different driving by initial and/or source data:

Equation (12)

Examples of observation operators are

for some points ${x}_{i}\in \bar{{\Omega}}$ and some weight functions ηi L(ω), Σi ⊆ ∂Ω, provided that these evaluations are well-defined, that is, $u(t)\in C(\bar{{\Omega}})$ in case (a), or u(t) ∈ L1(ω) in case (b), which according to theorem 2.1 is the case if ${\Omega}\subseteq {\mathbb{R}}^{d}$ with d = 1 for case (a) or d ∈ {1, 2, 3} for case (b).

The coefficients in (5) that we aim to in recover are

Moreover, we consider the problem of identifying, besides the differentiation orders, some space dependent quantities. In practice the most relevant ones are either the speed of sound c = c(x) in ultrasound tomography, or the initial data u0 = u0(x) while u1 = 0 (equivalently the source term f = f(x), while u0 = 0, u1 = 0) in photoacoustic tomography pat. For the latter purpose, we will consider observations not only at single points (a) or patches (b), but over a surface Σ.

2. The forward problem

2.1. Analysis of the forward problem

In this section we prove well-posedness of the forward problem (6). Just as in the inverse problem which is the main topic of this paper, we focus on the case of constant coefficients bk , dj , while the sound speed c = c(x) contained in the operator $\tilde{\mathcal{A}}$ may vary in space. We refer to chapter 7 of the author' upcoming book [22] for an analysis in the case when all coefficients are space dependent.

Our analysis will be based on energy estimates obtained by multiplying the pde (or actually its Galerkin semidiscretization with respect to space) with appropriate multipliers. For this purpose, we will make use of certain mapping properties of fractional derivative operators that can be quantified via upper and lower bounds.

We will use the following results that hold for $\theta \in \left[0,1\right)$.

  • From [10, lemma 2.3]; see also [35, theorem 1]: for any wHθ/2(0, t),
    Equation (13)
  • From [23, theorems 2.1, 2.2, 2.4 and proposition 2.1]: there exist constants $0< \underline{C}(\theta )< \bar{C}(\theta )$ such that for any wHθ (0, t) (with w(0) = 0 if $\theta \geqslant \frac{1}{2}$), the equivalence estimates
    Equation (14)
    hold. Likewise, for any wL2(0, t) we have ${\partial }_{t}^{\theta }w\in {H}^{-\theta }(0,t)$ and
    Equation (15)
    Moreover, (by a duality argument) there exists a constant C(2θ) > 0 such that for any wL2(0, t)
    Equation (16)
    Note that the latter two facts (15) and (16) hold true for θ = 1 as well.

We are now prepared to state and prove our main theorem on well-posedness of an initial value problem

Equation (17)

which clearly covers the models from section 1.1 in case of β1 = ⋯ = βK . Here $\tilde{\mathcal{A}}$ is a selfadjoint (with respect to some possibly weighted space ${L}_{w}^{2}({\Omega})$ with inner product ⟨⋅, ⋅⟩) operator with eigenvalues and eigenfunctions ${({\lambda }_{i},{\varphi }_{i})}_{i\in \mathbb{N}}$, and ${\dot{H}}^{1}({\Omega})=\left\{v\in {L}_{w}^{2}({\Omega}):{\Vert}{\tilde{\mathcal{A}}}^{1/2}v{{\Vert}}_{{L}_{w}^{2}({\Omega})}< \infty \right\}$. This includes the cases $\tilde{\mathcal{A}}=-{\triangle}$, $\tilde{\mathcal{A}}=-c{(x)}^{2}{\triangle}$, $\tilde{\mathcal{A}}=-c{(x)}^{2}\nabla \cdot (\frac{1}{\rho (x)}\nabla )$ relevant for acoustics, where in the latter two cases we use the inner product $\langle {v}_{1},{v}_{2}\rangle ={\int }_{{\Omega}}\frac{1}{{c}^{2}(x)}{v}_{1}(x){v}_{2}(x)\mathrm{d}x$ in which $\;\tilde{\phantom{\rule{0ex}{0ex}}\mathcal{A}}$ is selfadjoint.

The coefficients and orders will be assumed to satisfy

Equation (18)

and

Equation (19)

with k* to be specified below. The condition γJ αK has been shown to be necessary for thermodynamic consistency, see [1, lemma 3.1].

Note that the proof of theorem 2.1 also goes through without the restrictions γj ⩽ 1, αk ⩽ 1, but orders larger than one would require introducing further initial conditions, which would make the exposition less transparent. Moreover, the focus in this paper is actually on wave type equations where this restriction of the orders is physically relevant.

We first of all consider the following time integrated and therefore weaker form with $\tilde{r}(x,t)={\int }_{0}^{t}r(x,s)\mathrm{d}s$

Equation (20)

The results are obtained by testing the pde in (17) with ${\partial }_{t}^{\tau }u$ for some τ > 0 that needs to be well chosen in order to exploit the estimates (13)–(16). Generalizing the common choice τ = 1 in case of the classical second order wave equation γ0 = ⋯ = γJ = α0 = ⋯ = αK = 0 and in view of the stability condition γJ αK we choose τ ∈ [1 + γJ , 1 + αK ], actually, $\tau =1+{\alpha }_{{k}_{\ast }}$ with k* according to

Equation (21)

if α0 < γJ and k* = 0 otherwise. More precisely, we will assume that

Equation (22)

holds. Additionally, we assume

Equation (23)

The upper and lower bounds on ${\alpha }_{{k}_{\ast }}$ in (23) are needed to guarantee $\theta \in \left[0,1\right)$ in several instances of θ used in (13)–(16) in the proof below.

Note that we are considering the case γJ = αK only in the specific setting of a single $\tilde{\mathcal{A}}$ term. Indeed, when admitting several $\tilde{\mathcal{A}}$ terms, well-posedness depends on specific coefficient constellations. As an example of this, let us mention the fractional Moore–Gibson–Thompson equations arising in acoustics

Equation (24)

with positive parameters $\mathfrak{T},\mathfrak{b},\mathfrak{c},\alpha $. It is well-posed only if $\mathfrak{b}\geqslant {\mathfrak{c}}^{2}\mathfrak{T}$, according to [20, proposition 7.1]. As the analysis of this example also shows, even lower order time derivatives of $\tilde{\mathcal{A}}u$ (such as the $-{\mathfrak{c}}^{2}{\Delta}u$ term in (24)) may prohibit global in time boundedness of solutions. It will turn out in the proof that we therefore have to impose a bound on the magnitude of the lower order term coefficients in terms of the highest order one

Equation (25)

that due to T dependence of ${{\Vert}}_{0}{I}_{t}^{({\alpha }_{K}-{\alpha }_{k})/2}{{\Vert}}_{{L}^{2}(0,T)\to {L}^{2}(0,T)}$ clearly restricts the maximal time T unless all bk with index lower than k* vanish, cf corollary 2.1 below.

By applying the test functions as mentioned above, we will derive estimates for the energy functionals (which are not supposed to represent physical energies but are just used for mathematical purposes)

Equation (26)

where

Equation (27)

where $\tilde{d}={d}_{J}\,\mathrm{cos}(\pi (1+{\gamma }_{J}-{\alpha }_{{k}_{\ast }})/2)\underline{C}{((1+{\gamma }_{J}-{\alpha }_{{k}_{\ast }})/2)}^{2}$ with $\underline{C}$ as in (15), $\theta =(1+{\gamma }_{J}-{\alpha }_{{k}_{\ast }})/2\in \left[0,1\right)$ and

Equation (28)

where the sum in (28) is void in case k* = K.

The solution space induced by these energies is Uγ Uα where

Equation (29)

for any ɛ ∈ (0, αK − max{αK−1, γJ }).

Theorem 2.1. Assume that the coefficients γj , αk , dj , bk satisfy (18), (19), (22), (23), and that T and/or the coefficients bk for kk* are small enough so that (25) holds.

Then for any ${u}_{0},{u}_{1}\in {\dot{H}}^{1}({\Omega})$, (${u}_{2}\in {L}_{w}^{2}({\Omega})$ if γJ > 0), $r\in {H}^{{\alpha }_{{k}_{\ast }}-{\gamma }_{J}}(0,T;{L}_{w}^{2}({\Omega}))$, the time integrated initial boundary value problem (20) (to be understood in a weak ${\dot{H}}^{-1}({\Omega})$ sense with respect to space and an L2(0, T) sense with respect to time) has a unique solution uUγ Uα cf (29). This solution satisfies the energy estimate

Equation (30)

for some constant C(t) depending on time, which is bounded as t → 0 but in general grows exponentially as t, ${\mathcal{E}}_{\gamma }$, ${\mathcal{E}}_{\alpha }$ defined as in (26), (28), and

Equation (31)

Proof. We exclude the case αK = 0 since this via (18) implies that also γJ = 0 and we would therefore deal with the conventional second order wave equation whose analysis can be found in textbooks, cf e.g., [11]. Thus for the remainder of this proof we assume αK > 0 to hold.

Step 1. Galerkin discretisation.

In order to prove existence and uniqueness of solutions to (17), we apply the usual Faedo–Galerkin approach of discretisation in space with eigenfunctions φi corresponding to eigenvalues λi of $\tilde{\mathcal{A}}$, $u(x,t)\approx {u}^{L}(x,t)={\sum }_{i=1}^{L}\,{u}_{i}^{L}(t){\varphi }_{i}(x)$ and testing with φj , that is,

Equation (32)

To prove existence of a solution ${\underline{u}}^{L}\in {H}^{2+{\gamma }_{J}}(0,T;{\mathbb{R}}^{L})$ to (32), we rewrite it as a system of Volterra integral equations for

Unique solvability of the resulting system in L2(0, TL ) follows from [16, theorem 4.2, p 241 in section 9]. Then from

in case γJ > 0 we have a unique ${\underline{u}}_{tt}^{L}\in {H}^{{\gamma }_{J}}(0,{T}^{L})$ (cf [23, section 3.3]); the same trivially holds true in case γJ = 0.

From the energy estimates below, this solution actually exists for all times $t\in \left[0,T\right)$, so the maximal time horizon of existence is actually TL = T, independently of the discretisation level L.

Step 2. Energy estimates.

With k* defined by (21) in case (i) and set to k* = 0 in case (ii), we use $v={\partial }_{t}^{1+{\alpha }_{{k}_{\ast }}}{u}^{L}(s)$ as a test function in (32) (with t replaced by s) and integrate for s from 0 to t. The resulting terms can be estimated as follows.

Inequality (13) with $\theta =1+{\gamma }_{j}-{\alpha }_{{k}_{\ast }}\in \left[0,1\right)$ yields

and with $\theta =1+{\alpha }_{{k}_{\ast }}-{\alpha }_{k}\in \left[0,1\right)$

In particular for k = K > k* (which implies $1+{\alpha }_{{k}_{\ast }}-{\alpha }_{K}< 1$; the case k* = K will be considered below), by (15) with $\theta =(1+{\alpha }_{{k}_{\ast }}-{\alpha }_{K})/2\in \left[0,1\right)$, noting that ${\partial }_{t}^{(1+{\alpha }_{{k}_{\ast }}+{\alpha }_{K})/2}{\tilde{\mathcal{A}}}^{1/2}{u}^{L}$ vanishes at t = 0 and that ${\partial }_{t}^{(1+{\alpha }_{{k}_{\ast }}-{\alpha }_{k})/2}{\partial }_{t}^{(1+{\alpha }_{{k}_{\ast }}+{\alpha }_{K})/2}{\tilde{\mathcal{A}}}^{1/2}{u}^{L}={\partial }_{t}^{1+{\alpha }_{{k}_{\ast }}}{\tilde{\mathcal{A}}}^{1/2}{u}^{L}$,

Equation (33)

with

From the identity ${\int }_{0}^{t}\langle {\partial }_{t}v(s),v(s)\rangle \mathrm{d}s=\frac{1}{2}{\Vert}v(t){{\Vert}}^{2}-\frac{1}{2}{\Vert}v(0){{\Vert}}^{2}$ and Young's inequality in the case ${\gamma }_{J}={\alpha }_{{k}_{\ast }}$ we get

and

Hence in the case k* = K, where due to our assumption that αK > 0 at the beginning of the proof, ${\partial }_{t}^{{\alpha }_{K}}{\tilde{\mathcal{A}}}^{1/2}{u}^{L}(0)=0$,

Since the difference $1+{\alpha }_{{k}_{\ast }}-{\alpha }_{k}$ is larger than one for k < k*, we have to employ (16) in that case (for this purpose, note that still, according to our assumptions, $\theta =(1+{\alpha }_{{k}_{\ast }}-{\alpha }_{k})/2\in [0,1]$). From this and (15), again with $\theta =(1+{\alpha }_{{k}_{\ast }}-{\alpha }_{k})/2\in [0,1]$, we obtain

Equation (34)

Note that in case k* < K, the potentially negative contributions arising from the k < k* terms (34) are dominated by means of the leading energy term from (28), that is, ${\Vert}{\partial }_{t}^{(1+{\alpha }_{{k}_{\ast }}+{\alpha }_{K})/2}{\tilde{\mathcal{A}}}^{1/2}{u}^{L}{{\Vert}}_{{L}^{2}(0,t;{L}_{w}^{2}({\Omega}))}^{2}$ from (33), due to the fact that $(1+{\alpha }_{{k}_{\ast }}+{\alpha }_{k})/2< (1+{\alpha }_{{k}_{\ast }}+{\alpha }_{K})/2$ for k < k*.

To achieve this also in case k* = K with αK−1 < γJ < αK , we additionally test with $v={\partial }_{t}^{1+{\alpha }_{K}-\varepsilon }{u}^{L}(s)$ in (32) where 0 < ɛ < αK − max{αK−1, γJ }. Repeating the above estimates with ${\alpha }_{{k}_{\ast }}$ replaced by αK ɛ one sees that the leading energy term then is ${\Vert}{\partial }_{t}^{{\alpha }_{K}+(1-\varepsilon )/2}{\tilde{\mathcal{A}}}^{1/2}{u}^{L}{{\Vert}}_{{L}^{2}(0,t;{L}_{w}^{2}({\Omega}))}^{2}$ from (33). This still dominates the potentially negative contributions arising from the k < K terms according to (34) (with αK ɛ in place of ${\alpha }_{{k}_{\ast }}$), since (1 + αK ɛ + αk )/2 < αK + (1 − ɛ)/2 for k < K.

Note that in the case γJ = αK = α0, no potentially negative terms (34) arise.

Finally, the right-hand side term can be estimated by means of (15) with $\theta ={\alpha }_{{k}_{\ast }}-{\gamma }_{J}\in \left[0,1\right)$ and Young's inequality

with some constants $\bar{d},C > 0$.

Altogether we arrive at the following estimate

Equation (35)

where ${\mathcal{E}}_{\gamma }$, ${\underline{\mathcal{E}}}_{\gamma }$, ${\mathcal{E}}_{\alpha }$ are defined as in (26)–(28),

where the Boolean variable 1B is equal to one if B is true and vanishes otherwise.

We now substitute ${\mathcal{E}}_{\gamma }[u](t)$, ${\mathcal{E}}_{\alpha }[u](t)$ by lower bounds containing only an estimate from below of their leading order terms ${\underline{\mathcal{E}}}_{\gamma }[u](t)$ as defined in (27) and

Equation (36)

We then obtain from (35)

Equation (37)

where

It is readily checked that $\text{rhs}[u](t)\leqslant {c}_{b,\alpha }(t){\underline{\mathcal{E}}}_{\alpha }[u](t)$ for

Therefore, under the smallness assumption (25), that is, cb,α (T) < 1 estimate (37) yields

for some $\tilde{C}$ independent of t, L, with C0 as in (31). Applying Gronwall's inequality and re-inserting into (35) we end up with

Equation (38)

Step 3. Weak limits.

As a consequence of (38), the Galerkin solutions uL are uniformly bounded in Uγ Uα . Therefore the sequence ${({u}^{L})}_{L\in \mathbb{N}}$ has a weakly(*) convergent subsequence ${({u}^{{L}_{k}})}_{k\in \mathbb{N}}$ with limit uUγ Uα . To see that u satisfies (20), we revisit (32), integrate it with respect to time and take the L2(0, T) product with arbitrary smooth test functions ψ to conclude that

for all v ∈ span(φ1, ..., φK ), KLk and any $\psi \in {C}_{c}^{\infty }(0,T)$. Taking the limit k by using the previously mentioned weak(*) limit of ${u}^{{L}_{k}}$ we conclude

for all v ∈ span(φ1, ..., φK ) with arbitrary $K\in \mathbb{N}$, and any $\psi \in {C}_{c}^{\infty }(0,T)$. Therefore the weak limit u indeed satisfies the time integrated pde and we have proven the existence part of the theorem.

To verify the initial conditions, we use the fact that due to our assumption αK > 0, Uα continuously embeds into $C([0,T];\dot{H}({\Omega}))$ and therefore u(0) = u0 is attained in an $\dot{H}({\Omega})$ sense. Also, Uγ continuously embeds into ${C}^{1}([0,T];{L}_{w}^{2}({\Omega}))$ in case ${\gamma }_{J}< {\alpha }_{{k}_{\ast }}$ or γJ > 0. In the remaining case ${\gamma }_{J}={\alpha }_{{k}_{\ast }}=0$, attainment of ut (0) = u1 in an ${L}_{w}^{2}({\Omega})$ sense can be shown analogously to the conventional second order wave equation, cf, e.g., [11, theorem 3 in section 7.2].

Moreover, taking weak limits in the energy estimate (38) together with weak lower semicontinuity of the norms contained in the definitions of ${\mathcal{E}}_{\gamma }$, ${\mathcal{E}}_{\alpha }$, implies (30).

Step 4. Uniqueness.

The manipulations carried out in step 2 of the proof are also feasible with the Galerkin approximation uL replaced by a solution u of the pde itself, and lead to the energy estimate (30) independently of the Galerkin approximation procedure. From this we conclude that the pde with vanishing right-hand side and initial data only has the zero solution, which due to linearity of the pde implies uniqueness.□

Corollary 2.1. Under the assumptions of theorem 2.1 with bk = 0, k = 0, ..., k* − 1, the assertion extends to hold globally in time, that is, for all t ∈ (0, ). If additionally, r = 0, then the energy is globally bounded

with some constant C > 0 independent of time.

Remark 2.1. Using the so-called multinomial (or multivariate) Mittag–Leffler functions and separation of variables in principle enables a solution representation by separation of variables, cf [26, theorem 4.1]. However, proving convergence of the infinite sums in this representation would require extensive estimates of these Mittag–Leffler functions. Thus, in order to keep the exposition transparent without having to introduce too much additional machinery, we chose to remain with the energy argument in the proof above.

In order to establish sufficient regularity of the highest order term ${\partial }_{t}^{2+{\gamma }_{J}}{u}^{L}$ so that also the original equation (17) holds, we make use of the pde itself.

Corollary 2.2. Under the conditions of theorem 2.1 the solution u of (20) satisfies ${\partial }_{t}^{2+{\gamma }_{J}}u\in {L}^{2}(0,T;\dot{H}{({\Omega})}^{\ast })$ and is therefore also a solution to the original pde (17) in an ${L}^{2}(0,T;\dot{H}{({\Omega})}^{\ast })$ sense.

Proof. From the Galerkin approximation (32) of (17), due to invariance of the space span(φ1, ..., φL ) we can substitute v by ${\tilde{\mathcal{A}}}^{-1/2}v$ there and move ${\tilde{\mathcal{A}}}^{-1/2}$ to the left-hand side of the inner product by taking the adjoint to arrive at

Equation (39)

From theorem 2.1 we know that $\tilde{r}=-{\sum }_{k=0}^{K}\,{\partial }_{t}^{{\alpha }_{k}}{\tilde{\mathcal{A}}}^{-1/2}[{b}_{k}\tilde{\mathcal{A}}u(t)]+{\tilde{\mathcal{A}}}^{-1/2}r(t)$ is contained in ${L}^{2}(0,T;{L}_{w}^{2}({\Omega}))$. Thus (39) can be viewed as a Galerkin discretisation of a Volterra integral equation for $\zeta ={\partial }_{t}^{2+{\gamma }_{J}}{\tilde{\mathcal{A}}}^{-1/2}u(t)$ with $\tilde{r}$ as right-hand side. The corresponding coefficient vector ${\underline{\zeta }}^{L}$ of the solution ζL is therefore bounded by

with C independent of L. From this we deduce

and thus, by taking the limit L the same bound for ${\Vert}{\zeta }^{L}{{\Vert}}_{{L}^{2}(0,T;{L}_{w}^{2}({\Omega}))}$ itself.□

2.2. Numerical solution of the forward problem

There are several methods available for solving multi-term fractional ordinary differential equations. For our situation we require a solver that is accurate over the entire time interval $\left[0,\infty \right)$ since one problem we consider is the recovery of the components of the operator from large time values only. Our method of choice here is to convert the system with multiple fractional derivatives to a single ode of the form ${D}_{t}^{\gamma }u=Au$ where A is a square matrix. The seminal starting point here is the paper by Diethelm and Ford [7], which allows one to consider the equation ${D}^{{\alpha }_{0}}y(t)=f\left(t,y(t),{D}^{{\alpha }_{1}},\dots ,{D}^{{\alpha }_{n}}\right)$ subject to ${y}^{(k)}(0)={y}_{0}^{k}$ for k = 0, 1, ... ⌈α0⌉ − 1, and its conversion to the form ${D}^{\gamma }Y(t)=g\left(t,Y(t)\right)$, with Y(0) = Y0. Existence and uniqueness results are obtained and a numerical method for the solution formulated. This was later refined by Garrappa and Popolizio [15], in the case of linear systems into a solver whose accuracy approaches machine precision. This is based on using the Mittag–Leffler function with a matrix argument A that uses different approaches depending of the location of the eigenvalues of A in the complex plane and combined into an algorithm that works for matrices with any eigenvalue locations. See also [28]. We describe this above approach briefly since it comes with a certain restriction that should be clarified.

The following result is proven in [28]:

Theorem 2.2. Consider the initial value problem

Equation (40)

where α0 > α1 > ⋯ αk−1 and 0 < αk−1 < 1. Assume that each ${\alpha }_{j}\in \mathbb{Q}$ and let M be the least common denominator of α0, ..., αk−1. Set γ = 1/M and N = 0. Then the initial value problem (41) is equivalent to the system of equations

Equation (41)

together with the initial conditions

Equation (42)

As a result of the above theorem the linear system can be reformulated to the form

Equation (43)

where the coefficient matrix $A\in {\mathbb{R}}^{N\times N}$ is of companion type.

Then it is possible to express the solution to (43) in terms of Mittag–Leffler functions with matrix arguments as

Equation (44)

For general matrix functions the Schur–Parlett algorithm is the method of choice and to this end the final component needed is an accurate evaluation of the matrix Mittag–Leffler function and its derivatives. This is both nontrivial and quite technical and we simply refer to the above references [15, 28], and also to [14] for the details. Code to implement what is needed for this step is available from Roberto Garrappa's homepage [13].

Note that in our case we have α0 = 2 and also the requirement that each of the exponents αj = qj /pj must be rational numbers. The above restrictions can lead to a quite large system of size lcm(p1, ..., pk−1) and hence a small value of γ. However, we use this solver only for the purpose of providing simulated data and our inversion techniques are independent of its use.

Since we are using asymptotic formulas it is essential to use a reliable and accurate solver that is capable to cover a wide range of time values. The code from [14, 15, 28] provably fulfils these requirements and our numerical experiments confirmed this.

2.3. Resolvent equations

Separation of variables based on the eigensystem of $\mathcal{A}$ yields a solution representation

and convergence of the sums is guaranteed by the energy estimates from theorem 2.1. Here wfn , w2n , w1n , w0n denote the solutions of the resolvent equation

From the Laplace transformed resolvent equations we obtain the resolvent solutions

with

where $\hat{w}$ denotes the Laplace transform of w.

Formally we can write

and therefore

3. Uniqueness approaches

Choose, for simplicity the excitations u0i = 0, u1i = 0, u2i = 0, fi ≠ 0 and assume that for two models given by

Equation (45)

we have equality of the observations

from two possibly different space parts of the excitation fi , ${\tilde{f}}_{i}$, and possibly using two different test specimen characterised by $\mathcal{A}$, $\tilde{\mathcal{A}}$, while the temporal part σ of the excitation is assumed to be the same in both experiments, and its Laplace transform to vanish nowhere $\hat{\sigma }(s)\ne 0$ for all s. By analyticity (due to the multinomial Mittag–Leffler functions representation cf [26, theorem 4.1]) we get equality for all t > 0, thus equality of the Laplace transforms

Equation (46)

in some sector ${\mathbb{C}}_{\theta }$ of the complex plane, due to our assumption that $\hat{\sigma }(s)=\hat{\tilde{\sigma }}(s)$ and $\hat{\sigma }(s)\ne 0$ for all s. We have to disentangle the sum over n with the asymptotics for s in order to recover the unknown quantities in (45) or at least some of them.

3.1. Smooth excitations

We make use of the Weyl estimate λn Cd n2/d in ${\mathbb{R}}^{d}$ as well as additional regularity of the excitation to approximate the sum over n (46) of fractions by a single fraction of the form

for large real positive s. We focus exposition on the source excitation case u = 0, = 0, 1, 2, f ≠ 0 from (46), the other cases follow similar. Also, since we only need one excitation here I = 1, we just skip the subscript i. Indeed, for $s\in \mathbb{R}$, s ⩾ 1 we can estimate the difference

where we have used Young's inequality to estimate

Since λn n2/d , the sum converges for $2(\mu -\frac{1}{p}) > \frac{d}{2}$. On the other hand, to get an $O({s}^{-(2+{\gamma }_{J}+\varepsilon )})$ estimate for D with ɛ > 0, we need 2 − αN + 2/p > 2 + γJ and therefore choose

Equation (47)

(taking into account the tilde version of the above estimate as well). This yields

with $\varepsilon =\mu -\frac{d}{4}-\frac{\mathrm{max}\left\{{\alpha }_{N}+{\gamma }_{J},{\tilde{\alpha }}_{\tilde{N}}+{\tilde{\gamma }}_{\tilde{J}}\right\}}{2} > 0$. Thus, our assumptions on μ, ν are

Equation (48)

In case of rational powers

Equation (49)

setting $z={s}^{1/\bar{q}}$ we read (51) as an equality between two rational functions of z, whose coefficients and powers therefore need to coincide (provided $\left\{{s}^{1/\bar{q}}:s\in {\mathbb{C}}_{\theta }\right\}\supseteq [\underline{z},\infty )$ for some $\underline{z}\geqslant 0$). In the general case of real powers one can use the induction proof from [18, theorem 1.1] to obtain the same uniqueness. This yields

From this it becomes clear that we get uniqueness of

provided all ${\beta }_{k}={\tilde{\beta }}_{k}$ are known and ${\lambda }_{1}={\tilde{\lambda }}_{1}$ (but still possibly unknown) since we can then divide by λ1 and ${\lambda }_{1}^{{\beta }_{k}}$ in (II) and (IV), respectively. However, there seems to be no chance to simultaneously obtain bk and βk from (IV). Therefore, we will consider a setting with two different excitations in section 3.2 for this purpose.

The same derivation can be made for the variable c = c(x) model (6) in place of (5), with c, $\mathcal{A}$, L2(Ω), ${\Vert}\varphi {{\Vert}}_{{L}^{2}}=1$ replaced by 1, ${\mathcal{A}}_{c}$, ${L}_{1/{c}^{2}}^{2}({\Omega})$, ${\Vert}\varphi {{\Vert}}_{{L}_{1/{c}^{2}}^{2}}=1$. From (II) with c, $\tilde{c}$ replaced by unity, we get ${\lambda }_{1}={\tilde{\lambda }}_{1}$ so that we do not need to assume this.

Theorem 3.1. Let $f,\tilde{f}\in {\dot{H}}^{2\mu +2\nu }({\Omega})$ for μ, ν sufficiently large so that (48) holds.

Then

  • (a)  
    For the solutions u, $\tilde{u}$ of (5) with vanishing initial data and known equal ${\beta }_{k}={\tilde{\beta }}_{k}$, possibly unknown but equal ${\lambda }_{1}={\tilde{\lambda }}_{1}$, possibly unknown (rest of) $\mathcal{A}$, $\tilde{\mathcal{A}}$, and possibly unknown $f,\tilde{f}$ implies
  • (b)  
    For the solutions u, $\tilde{u}$ of (6) with vanishing initial data and known equal ${\beta }_{k}={\tilde{\beta }}_{k}$, and possibly unknown ${\mathcal{A}}_{c}$, ${\tilde{\mathcal{A}}}_{\tilde{c}}$ implies

Remark 3.1. This uniqueness approach also extends to different excitation combinations such as u0 ≠ 0, uj = 0, j = 1, 2, f = 0; u1 ≠ 0, uj = 0, j = 0, 2, f = 0; u2 ≠ 0, uj = 0, j = 0, 1, f = 0; (the latter only in case γJ > 0).

Consider for example the case u0 ≠ 0, uj = 0, j = 1, 2, f = 0 that will be relevant for additionally recovering u0, cf remark 3.6 and which is the one that differs most from the setting above in view of the fact that the numerator of the resolvent solution depends on s and λn . The crucial estimate on the error that is made by replacing λn by λ1 then becomes, for $s\in \mathbb{R}$, s ⩾ 1, and using the decomposition ω(s, λ) = s2 + ωd (s) + ωb (s, λ) with ${\omega }_{d}(s)={\sum }_{j=1}^{J}{\,d}_{j}\,{s}^{{\gamma }_{j}+2}$, ${\omega }_{b}(s,\lambda )={\sum }_{k=1}^{N}\,{s}^{{\alpha }_{k}}{b}_{k}{\lambda }^{{\beta }_{k}}$

for $\bar{d}={\sum }_{j=1}^{J}{\,d}_{j}$, $\bar{b}={\sum }_{k=1}^{N}\,{b}_{k}\,\mathrm{max}\left\{1,{\lambda }_{1}\right\}$, where we have used (λn λ1)(s2 + ωd (s)) ⩾ 0 and ${\lambda }_{n}{\omega }_{b}(s,{\lambda }_{1})-{\lambda }_{1}{\omega }_{b}(s,{\lambda }_{n})={\sum }_{k=1}^{N}\,{s}^{{\alpha }_{k}}{b}_{k}\left({\lambda }_{n}^{{\beta }_{k}}{\lambda }_{1}^{{\beta }_{k}}({\lambda }_{n}^{1-{\beta }_{k}}-{\lambda }_{1}^{1-{\beta }_{k}})\right)\geqslant 0$.

To achieve an $O({s}^{-(2+\mathrm{max}\left\{{\gamma }_{J},{\tilde{\gamma }}_{\tilde{J}}\right\}+\varepsilon )})$ estimate, instead of (48) we thus assume

Equation (50)

to recover theorem 3.1 with f, $\tilde{f}$ replaced by u0, ${\tilde{u}}_{0}$.

Remark 3.2. Note that in view of the fact that $\mathcal{A}$, ${\mathcal{A}}_{c}$ do not need to be known (up to the fact that a single eigenvalue is supposed to coincide), we are able to identify all relevant information on the damping model in an unknown medium.

3.2. Single mode excitation

Assume that we know at least some of the eigenfunctions and can use them as excitations ${f}_{i}={\varphi }_{{n}_{i}}$, ${\tilde{f}}_{i}={\tilde{\varphi }}_{{\tilde{n}}_{i}}$, i = 1, ..., I, so that (46) becomes

Equation (51)

We can use the rational function or induction argument as above to compare powers of s and conclude the following:

To extract bk and βk , which can be done separately for each k, (skipping the subscript k) we use I = 2 excitations, assume that the corresponding eigenvalues of $\mathcal{A}$ and $\tilde{\mathcal{A}}$ are known and equal and define

One easily sees that the 2 × 2 matrix F'(b, β) is regular for ${\lambda }_{{n}_{1}}\ne {\lambda }_{{n}_{2}}$ and b ≠ 0 and thus by the inverse function theorem F is injective. Thus from (I)–(IV) we obtain all the constants in (45).

Theorem 3.2. Let ${f}_{1}={\varphi }_{{n}_{1}}$, ${f}_{2}={\varphi }_{{n}_{2}}$, ${\tilde{f}}_{1}={\tilde{\varphi }}_{{\tilde{n}}_{1}}$, ${\tilde{f}}_{2}={\tilde{\varphi }}_{{\tilde{n}}_{2}}$ for some n1n2, ${\tilde{n}}_{1}\ne {\tilde{n}}_{2}\in \mathbb{N}$. Then

for the solutions ui , ${\tilde{u}}_{i}$ of (5) with vanishing initial data, known and equal ${\lambda }_{i}={\tilde{\lambda }}_{i}$, i = 1, 2 and possibly unknown (rest of) $\mathcal{A}$, $\tilde{\mathcal{A}}$ implies

Remark 3.3. Scaling with unity in the excitation has been chosen here just for simplicity of exposition. Clearly also multiples of eigenfunctions can be used.

Again, uniqueness can be shown analogously with eigenfunction excitation by initial data rather than by sources.

We briefly consider the additional recovery of a spatially varying wave speed c and for this purpose look at the ch model (9) as an example. With ${\mathcal{A}}_{c}=-c{(x)}^{2}{\triangle}$, (I)–(IV) (with N = 1, c = 1, β = 1, bk = b) yields

Using inverse Sturm–Liouville theory, one can recover the spatially varying coefficient c(x) in one space dimension from knowledge of all eigenvalues (actually just two sets of eigenvalues corresponding to different boundary conditions), see, [6, 29]. However, in order to do so via this single mode excitation approach, we would need infinitely many measurements $\left\{{n}_{i}:i=1,\dots ,I\right\}=\mathbb{N}$.

The goal of obtaining knowledge on the eigenvalues can be much better achieved by going back to (46) and considering equality of poles (and residues), which according to [21, section 4] gives us λn and $b{\lambda }_{n}^{\beta }$ (note that we have set c = 1 since c(x) is contained in ${\mathcal{A}}_{c}$) from a single excitation u1 with nonvanishing Fourier coefficients. We will now consider this approach for the more general setting (3) in section 3.3.

3.3. Poles and residues

In this section we focus on the second order in time case (3) and achieve separation of summands in the eigenfunction expansion by using poles (and residues) of $\hat{h}(s)$. For this purpose we have to prove that each summand λn gives rise to at least one pole (in case of ch we know there are two complex conjugates) and that these poles differ for different λn . This is also known for ch and fz, see [21, section 4]. Indeed the results from [21, lemma 4.1, remark 4.1] on ch carry over to the model (3) in the setting

Equation (52)

that is, (3) with a single β and dissipative behaviour (as is natural to demand from a damping model). Existence of at least one pole for each λn can even be shown for the most general model (7).

Lemma 3.1. For each λn there exists at least one root of ${\omega }_{n}^{\text{mg}}(s)={s}^{2}+{c}^{2}{\lambda }_{n}+{\sum }_{k=1}^{N}{s}^{{\alpha }_{k}}{\sum }_{\ell =1}^{{M}_{k}}{c}_{k\ell }{\lambda }_{n}^{{\beta }_{k\ell }}$, that is, a pole of the Laplace transformed relaxation solution ${\hat{w}}_{fn}^{\text{mg}}$, ${\hat{w}}_{0n}^{\text{mg}}$, or ${\hat{w}}_{1n}^{\text{mg}}$. Moreover, the poles are single.

In case of (52), the poles of ${\hat{w}}_{fn}^{\text{g}}$, ${\hat{w}}_{0n}^{\text{g}}$, ${\hat{w}}_{1n}^{\text{g}}$, which are the roots of ${\omega }_{n}^{\text{g}}(s)={s}^{2}+{c}^{2}{\lambda }_{n}+{\sum }_{k=1}^{N}{b}_{k}{s}^{{\alpha }_{k}}{\lambda }_{n}^{\beta }$, lie in the left half plane and differ for different λn .

Proof. Let fn (z) = z2 + c2 λn , ${g}_{n}(z)={\sum }_{k=1}^{N}{z}^{{\alpha }_{k}}{\sum }_{\ell =1}^{{M}_{k}}{c}_{k\ell }{\lambda }_{n}^{{\beta }_{k\ell }}$. Let CR be the circle radius R, centre at the origin. Then, due to αk ⩽ 1, for a sufficiently large R > c2 λn , the estimate |gn (z)| < |fn (z)| holds on CR and so Rouché's theorem shows that fn (z) and ${\omega }_{n}^{\text{mg}}(z)=({f}_{n}+{g}_{n})(z)$ have the same number of roots, counted with multiplicity, within CR . For fn these are only at $z=\pm i\sqrt{{\lambda }_{n}}c$ and so ${\omega }_{n}^{\text{mg}}$ has precisely one single root in the third and in the fourth quadrant, respectively.

The fact that the poles of ${\hat{w}}_{fn}^{\text{g}}$, ${\hat{w}}_{0n}^{\text{g}}$, ${\hat{w}}_{1n}^{\text{g}}$ lie in the left half plane in case of (52) follows by a standard energy argument. Suppose now that ωg has a root at p = r eiθ , where π/2 ⩽ θ < π, for both λn and λm . Then for p = r eiθ , subtracting the equations ${\omega }_{n}^{\text{g}}(p)=0$ and ${\omega }_{m}^{\text{g}}(p)=0$ we obtain

Equation (53)

thus

Now if λn λm then the right-hand side is positive and real and so 2θ = π. This means that ${\sum }_{k=1}^{N}{b}_{k}{s}^{{\alpha }_{k}}={\sum }_{k=1}^{N}{b}_{k}{r}^{{\alpha }_{k}\pi /2}$ (where all bk have the same sign) has nonvanishing imaginary part, a contradiction to the fact that the left-hand side of (53) is real.□

Since we will separate the summands by taking residues on both sides of the identity (46), the following identity is also crucial.

Lemma 3.2. For each λn , $n\in \mathbb{N}$, the residues of ${\hat{w}}_{fn}^{\text{g}}$, ${\hat{w}}_{0n}^{\text{g}}$, ${\hat{w}}_{1n}^{\text{g}}$ do not vanish.

Proof. By l'Hospital's rule and due to the fact that the poles pn are single, we get for $\omega (s)={\omega }_{fn}^{\text{mg}}(s)={\omega }_{1n}^{\text{mg}}(s)$

Moreover, with $\omega (s)={\omega }_{0n}^{\text{mg}}(s)$, $d(s)=\frac{\omega (s)-c{\lambda }^{2}}{s}$,

As a consequence of lemma 3.1 and 3.2, by taking residues at the poles pn of the Laplace transformed time trace data $\hat{h}(s)$, we obtain from (46)

Equation (54)

where the cardinality of the index set ${K}^{{\lambda }_{n}}$ is equal to the multiplicity of the eigenvalue λn .

In order to achieve nonzero numerators in (46), (51) and (54), we assume

Equation (55)

and

Equation (56)

see remark 3.4 below.

Assume that,

Equation (57)

(e.g. by assuming $f=\tilde{f}$, $B=\tilde{B}$, $\mathcal{A}=\tilde{\mathcal{A}}$). Then from the reciprocals of (54) we get, after division by ${\sum }_{k\in {K}^{{\lambda }_{n}}}\langle f,{\varphi }_{n,k}\rangle B{\varphi }_{n,k}$ and subtracting 2pn on both sides

Equation (58)

In case $\beta ={\beta }_{k}={\tilde{\beta }}_{k}=\tilde{\beta }$, ${\lambda }_{n}={\tilde{\lambda }}_{n}$, this can be divided by ${({\lambda }_{n})}^{\beta }$ and simplifies to

that is, in case of rational powers (49) an equality between two polynomials of $z={s}^{1/\bar{q}}$ at infinitely many different (according to the second part of lemma 3.1) points ${p}_{n}^{1/\bar{q}}$. This yields

Thus we recover the results of theorem 3.1 under essentially more restrictive assumptions—however we can drop the smoothness assumption on the excitation f, so that f = δ is admissible. Thus the use of poles and residues has no clear advantage over the smooth excitation approach as long as we only aim at recovering the damping term constants. However, it allows to gain additional information on spatially varying coefficients c = c(x) or f = f(x) or u0 = u0(x).

3.3.1. Additional recovery of c(x)

First of all, just using knowledge of the pole locations of the data, lemma 3.1 allows to obtain the following corollary of theorem 3.1(b) on uniqueness of the numbers N, J, coefficients bk , orders αk , as well as the function c(x) in one space dimension. Note that the assumption f = δ is not compatible with the smoothness assumptions from theorem 3.1, thus we cannot apply the strategy from [21] of obtaining, besides the eigenvalues, also the values φ(x0) and therewith the norming constants of the eigenfunctions. We thus apply a different result from Sturm–Liouville theory, according to which a pair of eigenfunction sequences corresponding to two different boundary impedances suffices to recover the potential q of the differential operator −△ + q, see [6, 30]. Using the transformation of variables

so that

this transfers to the recovery of the sound speed c(x).

Since these results hold in one space dimension only, in this subsection we consider

Equation (59)

for some L > 0 and some nonnegative impedance constants H0, HL . For two different right hand boundary impedance values HL,i , i = 1, 2 we denote by ui the solution of (59) with HL = HL,i and by ${({\varphi }_{n,i},{\lambda }_{n,i})}_{n\in \mathbb{N}}$ the eigenpairs of −c(x)△ with boundary conditions ${\varphi }_{n,i}^{\prime }(0)-{H}_{0}{\varphi }_{n,i}(0)=0$, ${\varphi }_{n,i}^{\prime }(L)+{H}_{L,i}{\varphi }_{n,i}(L)=0$, and weighted L2 normalization ${\int }_{0}^{L}\frac{1}{c{(x)}^{2}}{\varphi }_{n,i}{(x)}^{2}\mathrm{d}x=1$, i = 1, 2.

For constructing sufficiently smooth excitations with nonvanishing coefficients with respect to the eigenfunction basis, we point to remark 3.4 below.

Moreover, we assume that the observations satisfy

Equation (60)

which is, e.g., the case with the choice Bi v := v(x0), x0 ∈ {0, L}, see remark 3.4 below.

Analogous notation is used for (59) with c, bk , αk , βk , f, H0, HL replaced by their corresponding tilde counterparts.

From theorem 3.1 and uniqueness of the eigenvalues from the poles we obtain the following result on combined identification of the coefficients in the damping model and the space-dependent sound speed.

Corollary 3.1. Let $f,\tilde{f}\in {\dot{H}}^{2\mu +2\nu }(0,L)$ for μ, ν satisfying (48) and let (56) and (60) hold.

Then

for the solutions ui , ${\tilde{u}}_{i}$ of (59) with vanishing initial data, HL = HL,i , ${\tilde{H}}_{L}={\tilde{H}}_{L,i}$, known equal ${\beta }_{k}={\tilde{\beta }}_{k}$, and unknown c, $\tilde{c}\in {C}^{2}(0,L)$, implies

The same result can be obtained without smoothness assumption on f provided (57), and $\beta ={\beta }_{k}={\tilde{\beta }}_{k}=\tilde{\beta }$ by a poles/residual argument, cf (58).

Remark 3.4. Some comments on the verification of conditions (56), (60) and (57) are in order, in particular in view of the fact that c(x), $\tilde{c}(x)$ and therefore also the eigenfunctions are unknown. First of all, note that choosing $f,\tilde{f}={\delta }_{{x}_{0}}$ formally satisfies this condition provided we can make sure that none of the eigenfunctions vanish at the point x0. This is, e.g., the case if x0 ∈ {0, L} is a boundary point for then in case x0 = L the impedance conditions would imply ${\varphi }_{n}^{\prime }(0)-{H}_{0}{\varphi }_{n}(0)=0$, ${\varphi }_{n}^{\prime }(L)=0$, φn (L) = 0, which by uniqueness of solutions to the second order ODE $-c(x){\varphi }_{n}^{\prime \prime }(x)-{\lambda }_{n}{\varphi }_{n}(x)=0$ with these boundary conditions would imply φn ≡ 0, a contradiction to φn being an eigenfunction. Likewise for x0 = 0.

Note that setting Bi v := v(x0), x0 ∈ {0, L}, the same argument yields (60) as well as (57).

However, clearly $f={\delta }_{{x}_{0}}$ does not satisfy the regularity requirements from the first part of the corollary, thus we also indicate a possibility of constructing a smooth f that satisfies (56). For any s > 0, ɛ > 0, the function $f(x)=c{(x)}^{2}{\sum }_{n=1}^{\infty }{\lambda }_{n}^{-s}{n}^{-(\frac{1}{2}+\varepsilon )}{\varphi }_{n}(x)$ lies in ${\dot{H}}^{s}$ and satisfies ${\langle f,{\varphi }_{k}\rangle }_{{L}_{1/{c}^{2}}^{2}}\geqslant {\lambda }_{n}^{-s}{n}^{-(\frac{1}{2}+\varepsilon )}$ for all $k\in \mathbb{N}$. Unfortunately this is only a result on existence of f, since the reconstruction requires knowledge of c.

3.3.2. Additional recovery of f(x) (or u0(x) or u1(x))

In order to reconstruct the possibly unknown excitation f (analogously for u1, u0, the latter relevant in PAT) we use measurements not only at finitely many points but on a surface Σ and make the linear independence assumption

Equation (61)

Moreover, for the purpose of recovering f (or u1, or u0) we assume c(x) to be known and therefore ${\varphi }_{n}={\tilde{\varphi }}_{n}$, ${\lambda }_{n}={\tilde{\lambda }}_{n}$. Thus in place of (54) we have

Equation (62)

Again, we can obtain a combined identification result from theorem 3.1, uniqueness of the eigenvalues from the poles, and (62).

Corollary 3.2. Let $f,\tilde{f}\in {\dot{H}}^{2\mu +2\nu }({\Omega})$ for μ, ν satisfying (48) and let (62) hold.

Then

for the solutions u, $\tilde{u}$ of (6) with vanishing initial data known equal ${\beta }_{k}={\tilde{\beta }}_{k}$, and unknown f, $\tilde{f}\in {L}^{2}({\Omega})$, implies

Remark 3.5. The same result can be obtained without smoothness assumption on f provided (57), and $\beta ={\beta }_{k}={\tilde{\beta }}_{k}=\tilde{\beta }$ by a poles/residual argument, cf (58).

Remark 3.6. The result carries over almost without change to the setting f = 0, u0 = 0, where u1 = u1(x) ≠ 0 is the unknown to be recovered.

In the setting f = 0, u1 = 0, of recovering u0 = u0(x) ≠ 0 (as relevant in PAT), we can rely on remark 3.1 to conclude that Corollary 3.2 remains valid with f, $\tilde{f}$, (48) replaced by u0, ${\tilde{u}}_{0}$, (50).

4. Reconstruction approaches

Throughout this section we will a single mode excitation and therefore skip the subscript ni in ${\varphi }_{{n}_{i}}$, ${\lambda }_{{n}_{i}}$ and i in Bi to simplify notation. Moreover, we focus on the second order model (3).

4.1. Full time observations or analyticity: Laplace domain reconstruction

Starting from the equation

(cf (51) for the slightly simpler setting of (5)) and noting that taking the reciprocal and rearranging makes the problem somewhat less nonlinear and the derivative easier to compute, we consider

with

and

For applying Newton's method, the Jacobian is given by

Taking the logarithm makes the equation less nonlinear with respect to α and β by considering (in case N = 1)

Then F1 is linear with respect to b and the Fm are linear with respect to α and β. This helps to enlarge the convergence radius of Newton's method, as we experienced in our numerical tests.

4.2. Large time asymptotics

In case u0 = φ, u1 = 0, f = 0 we have

Equation (63)

We will heavily make use of a Tauberian theorem (for s → 0, t, see [12, theorem 2 in chapter XIII.5]) which we here quote for the convenience of the reader.

Theorem 4.1. Let $w:\left[0,\infty \right)\to \mathbb{R}$ be a monotone function and assume its Laplace–Stieltjes transform $\hat{w}(s)={\int }_{0}^{\infty }{\mathrm{e}}^{-st}\,\mathrm{d}w(t)$ exists for all s in the right hand complex plane. Then for ρ ⩾ 0

if and only if

Note that here dw is actually a measure, which in case of an absolutely continuous function w can be written as dw(t) = w'(t)dt.

Applying theorem 4.1 to (63) we get

Equation (64)

Therefore we get an asymptotic formula (like the one in [17] for the subdiffusion equation) for the smallest order

Equation (65)

By l'Hospital's rule and actually thereby removing the constant we can instead also use

After having determined α1 this way, we can also compute b1 as a limit

Equation (66)

So we can successively (by the above procedure and subtracting the recovered terms one after another) reconstruct those terms k for which αk < 2α1, see the remainder term in (64). However, if there are terms with αk ⩾ 2α1, they seem to get masked by the $O({t}^{-2{\alpha }_{1}})$ remainder.

The same holds true if we do the excitation by u0 = 0, u1 = φ, f = 0, or u0 = 0, u1 = 0, f = φ, where (63) changes to $\hat{h}(s)=\frac{1}{{s}^{2}+{\sum }_{k=1}^{N}{b}_{k}{\lambda }^{{\beta }_{k}}{s}^{{\alpha }_{k}}+{c}^{2}\lambda }B\varphi $.

In order to avoid the restriction αk < 2α1, we thus refine the expansion of $\hat{h}(s)$ in terms of powers of s and retain the singular ones, that is, those with negative powers, since we are looking at the limiting case s → 0. Using the geometric series formula and the multinomial theorem, with the abbreviations ${\Sigma}={\sum }_{k=1}^{N}{b}_{k}{\lambda }^{{\beta }_{k}}{s}^{{\alpha }_{k}}$, $q=\frac{{s}^{2}+{\Sigma}}{{c}^{2}\lambda }$, yields

The terms corresponding to < m are obviously O(s), so to extract the singularities it suffices to consider = m. The set of indices leading to singular terms is then

and the singular part of $\hat{h}$ can be written as

Equation (67)

In case N = 2 with i = i1, i2 = mi, this reads as

since

For (67), by theorem 4.1

which confirms (65) and (66).

For the first damping term, we also get an asymptotic formula for the smallest order in Laplace domain—the small s counterpart to (65): due to

Equation (68)

we have

Equation (69)

One might realise this limit by extrapolation: fitting the values $\frac{\mathrm{log}\,\hat{h}({s}_{m})}{\mathrm{log}\,{s}_{m}}$ at M sample points s1, ..., sM to a regression line or low order polynomial r(s) and setting α1 = 1 + r(0). Next, recover b1 from

The formula (68) also shows why recovering b1 and α1 simultaneously appears to be so hard: the factor multiplied with log b1 is unity while the factor multiplied with 1 − α1 tends to infinity.

4.3. Small time asymptotics

In case u0 = φ, u1 = 0, f = 0 where the Laplace transformed observations are given by

Equation (70)

large s/small t asymptotics cannot be exploited for the following reason. As s then $\hat{h}(s)\to 0$ and in fact $s\hat{h}(s)\to 1$. More precisely, for

Thus clearly ${s}^{\gamma }[s\hat{h}(s)-1]\to 0$ for any $\gamma \in \left[0,2\right)$.

Not even large λ values (or some combined asymptotics as s → 0, λj ) seems to work because then the c2 λ term takes over and again masks the terms containing αk .

However, in case u0 = 0, u1 = φ, f = 0 the situation is different since then

Equation (71)

and therefore

Equation (72)

which via theorem 4.1 (where due to [12, theorem 3 in chapter XIII.5] we can replace s → 0, t by s, t → 0) yields

where we have used the fact that h(0) = 0, h'(0) = because of u0 = 0, u1 = φ.

This yields the asymptotic formulas

where we have used l'Hospital's rule and

Again, in order to also recover αN−1, ..., α1, one may also take into account mixed terms in the expansion of (72) analogously to the large t asymptotics case. In case of two damping terms, this results in

Equation (73)

5. Numerical reconstructions

In this section we give some details and show some reconstructions of the algorithms described in section 4. We assume single mode excitation, set βi = 1 and break the reconstructions into three separate paradigms as previously indicated depending on which time ranges are feasible to be measured.

5.1. Full time measurements

The leftmost graph in figure 1 shows the actual time trace of h(t) over the range 0 ⩽ t ⩽ 40 for a damping model with three terms ${b}_{i}{\partial }_{t}^{{\alpha }_{i}}$, i = 1, 2, 3. Note the damped oscillatory behaviour evident and the range shown is the one we are labelling as 'full-time' despite the fact that we will later look at the important case of the very long term behaviour of the solution where the time values are substantially outside of this range. This time range terminology is highly dependent on the values of critical constants in the model. These include the size of the domain (which determines the scaling of the eigenvalues {λn }), the wave speed c and of course the strength of the damping given by the coefficients {bi }. In our reconstructions we will typically take these constants to be of approximate order unity but note the dependence on these quantities and the effect that a substantial change would make.

Figure 1.

Figure 1. Profiles of h(t) and $\hat{h}(s)$, actual and after iterations 1 and 2.

Standard image High-resolution image

To obtain this h(t) our damping constants are of order $\sim 0.1-0.2$ and the combined term Λ = c2 λ is taken to be unity. Changes to the former would modify the time-decrease in h and changes to the latter would alter the frequency of the oscillations.

The rightmost graph in figure 1 shows the logarithm of the Laplace transform $H(s)=\hat{h}(s)$ together with the values of ${\hat{h}}_{i}(s)$ after iteration i of the scheme.

This demonstrates two important facts. First, the very fast convergence of the scheme in the sense of the convergence of the target Laplace transform function as the parameters {αi , bi } are resolved. The actual $\hat{h}(s)$ is shown by the solid black line, the initial approximation by the dotted red curve and the first iteration by the dashed green curve which, at this scale is already almost indistinguishable from the actual $\hat{h}$. Second, given this it should be quite feasible to obtain reasonably accurate values the parameters {αi , bi }. In addition, assuming we excite the system with an initial condition equal to an eigenfunction corresponding to an actual eigenvalue λ, we will be able to reconstruct the value of c2 from the composite Λ = c2 λ.

In this situation, to reconstruct the parameters {αi , bi , c2} we work directly with the representation $\hat{h}(s)$ or more exactly with its logarithm which is a more convenient form for computing the Jacobian

Equation (74)

Of course, to obtain $\hat{h}(s)$ we must approximate the integral $\hat{h}(s)={\int }_{0}^{\infty }{\mathrm{e}}^{-st}\,h(t)\mathrm{d}t$ and this indeed does require full time measurements. However, we do not actually need the values of the analytic function $\hat{h}(s)$ for all s for the Newton scheme and so if we only use s values with s > s0 then this allows us to ignore very large time measurements and indeed we truncated these to t ⩽ 40.

The reconstruction of the components in the function $\hat{h}(s)$ is shown graphically in figure 2 and in tabular form in table 1. The exact values were α = {0.25, 0.5, 0.75}, b = {0.2, 0.25, 0.1}, Λ = 4. Note that we have included the value of the residual here and keeping track of this value allows the iteration scheme to terminate when saturation has occurred.

Figure 2.

Figure 2. Reconstructed values of {bi } and {αi }. The symbols in red are the exact values.

Standard image High-resolution image

Table 1. Recovery of damping terms and unknown Λ from full time values.

Iter α1 α2 α3 b1 b2 b3 ΛResidual
00.30000.60000.80000.30000.37500.15003.500 
10.24480.55900.79460.19070.26120.10984.0030.035 032
20.25060.52540.76950.20570.27470.11904.0000.004 371
30.24920.52530.77000.20600.27840.11904.0000.000 075
40.24910.52540.77000.20600.27480.11924.0000.000 000

5.2. Large time measurements

Here we are trying to simulate the asymptotic values of the constituent powers of t occurring in the data function h(t). This was achieved by using a sample of points between tmin and tmax.

A Newton based approach

We apply Newton's method to recover the constants {pi , ck, } in

which results from the expansion (67) in case of three terms, and then recover αi = pi and bi = c1,i Γ(1 − pi ). In the above we have neglected the term ${t}^{-3{\alpha }_{3}}$ as this would not arise from the Tauberian theorem in the case that the largest power ${\alpha }_{3}\geqslant \frac{1}{3}$. We may also have to exclude other terms such as ${t}^{-2{\alpha }_{1}-{\alpha }_{2}}$ if 2α1 + α2 ⩾ 1. In practice, during the iteration process, terms should be included or excluded in the code depending on this criterion: we did so by checking if the argument passed to the Γ function would be negative in which case the term is deleted from use for that iteration step.

Values for the table shown below were tmin = 5 × 104 and tmax = 2 × 105. As a general rule, terms with small α values can be resolved with a smaller value of tmax, but for, say the recovery of a pair of damping terms with αi > 0.8, a larger value of tmax with commensurate accuracy will be needed.

The case of three damping terms ${\left\{{b}_{i}{\partial }_{t}^{{\alpha }_{i}}\right\}}_{i=1}^{3}$ with $\alpha =\left\{\frac{1}{4},\frac{1}{3},\frac{2}{3}\right\}$ and bi = 0.1 is shown in table 2 below. The initial starting values were taken to be between ten and thirty percent of the actual. These are shown in the line corresponding to iteration 0.

Table 2. Large time values with three damping terms.

Iter α1 α2 α3 b1 b2 b3 Residual
00.20000.30000.60000.1300.0800.110 
10.24090.35720.61680.1120.0890.1070.548 761
20.24770.33070.65310.1020.0910.1040.104 843
30.24990.33460.66500.1000.0990.0920.005 124
40.25000.33320.66680.1000.0990.0900.003 625
50.25000.33310.66580.1000.1000.0900.001 134
60.25000.33320.66580.1000.1000.0920.000 578
70.25000.33320.66600.1000.1000.0930.000 357
80.25000.33330.66650.1000.1000.0940.000 240

There are features here that are typical of such reconstructions. The reconstruction method resolves the lowest fractional power α1 and its coefficient b1 quickly as this term is the most persistent one for large times: essential numerical convergence for {α1, b1} is obtained by the third iteration. The next lowest power and coefficient lags behind; here {α2, b2} is already at the stated accuracy by the fifth iteration. In each case the power is resolved faster and more accurately than its coefficient. The third term also illustrates this; the power is essentially resolved by iteration 8, but in fact its coefficient b3 is not resolved to the third decimal place until iteration number 30. This is seen quite clearly in the singular values of the Jacobian: the largest singular values correspond to the lowest α-values and the smallest to the coefficients of the largest α-powers.

As might be expected, resolving terms whose powers are quite close is in general more difficult. This is relatively insignificant for low α values. For example, with α1 = 0.2 and 0.22 < α2 < 0.25 say, correct resolution will be obtained although the coefficients will take longer to resolve than indicated in table 2. On the other hand if, say ${\alpha }_{1}=\frac{1}{4}$ and α2 = 0.85, α3 = 0.9, then with the indicated range of time values used the code will fail to recover this last pair. If this is sensed and now only a single second power is requested this will give a good estimate for α2 but its coefficient will be overestimated. Also, as indicated previously, α values close to one require an extended time measurement range to stay closer to the asymptotic regime of u(t).

Note that the starting values were at least ten percent away from the exact values and one cannot expect the convergence radius of Newton's methods to be much larger for a problem exhibiting as high nonlinearity as the one at hand. Still, let us point to the fact that in the single term case we do not need any starting guesses but obtain very good results directly from the asymptotic formulas (65) and (66). This might give the idea of using the asymptotics to construct starting guesses in the multi term case and then apply Newton. However, it is unclear how the asymptotics could yield starting values for terms other than the first one. A (theoretical, as it turns out) possibility for exploiting asymptotics in the multi term case is described below.

Successive sequential use of asymptotic formulas

A few words are in order about an approach that from the above discussion might seem a good or even a better alternative.

Since each damping term ${b}_{i}{\partial }_{t}^{{\alpha }_{i}}$ contributes a time trace term with large time behaviour ${c}_{i}{t}^{-{\alpha }_{i}}$, it is feasible to take T sufficiently large so that ${c}_{i}{t}^{-{\alpha }_{i}}\ll {c}_{1}{t}^{-{\alpha }_{1}}$ for t > T, that is, all but the smallest damping power is negligible, and this can then be recovered. In successive steps then we subtract this from the data h to get ${h}_{1}(t)=h(t)-{c}_{1}{t}^{-{\alpha }_{1}}$ and now seek to recover the next lowest α power from the large time values of h1(t) in a range (δT, T) for 0 < δ < 1. Then these steps can be repeated until there is no discernible signal remaining in the sample interval tmin, tmax.

This indeed works well under the right circumstances for recovering two α values but the coefficients {bi } are less well resolved. It also requires a delicate splitting of the time interval and gives a much poorer resolution of the two terms in the case where, say α1 = 0.2 and α2 = 0.25 than that recovered from the Newton scheme. For the recovery of three damping terms this was in general quite ineffective.

Every time an αi has been recovered, the remaining signal is significantly smaller than the previous, leading to an equally significant drop in effective accuracy. Also, even if we just make a small error in the coefficient bi , the relative error that is caused by this becomes completely dominant for large times.

In short, this is an elegant and seemingly constructive approach to showing uniqueness for a finite number of damping terms. However, it has limited value from a numerical recovery perspective when used under a wide range of parameter values.

5.3. Small time measurements

In this case we are simulating measurements taken over a very limited initial time range: in fact we take the measurement interval to be $t\in \left[0,0.1\right)$. The line of attack is to use the known form of $\hat{h}(s)$ for large values of s and convert the powers of s appearing into powers of t for small times using the Tauberian theorem. In case of two damping terms this gives

where each term ck, is computed in terms of {αi }, {bi } and λ, cf (73).

The values of {αi , ck, } were then computed from the data by a Newton scheme then finally converted back to the derived values of {αi } and {bi }.

The exact values chosen were α = {0.25, 0.2}, b = {0.1, 0.1} and the initial starting guesses were α = {0.3, 0.16}, b = {0.08, 0.12}. We show the progression of the iteration in table 3.

Table 3. Small time values with two damping terms.

Iter α1 α2 b1 b2 Residual
00.30000.16000.08000.1200 
10.28480.16730.08210.11600.034 308
20.26320.19440.08630.11420.024 472
30.25540.20160.08620.11400.002 015
40.25370.20330.08610.11390.000 444
50.25340.20360.08610.11390.000 059
60.25340.20360.08610.11390.000 016
70.25340.20360.08610.11390.000 010

While theory predicts reconstructibility of an arbitrary number of terms in both cases, there is a clear difference in ability to reconstruct terms between the small time and the large time asymptotics. First of all, the method we described effectively only recovers two terms with small time measurements, as compared to three in the large time. The bi coefficients, which are always harder to obtain than the αi exponents, are much worse than in the small time than in the large time regime. This is partly explained by the higher degree of ill-posedness due to the necessity of differentiating the data twice.

Acknowledgments

The work of the first author was supported by the Austrian Science Fund fwf under the Grants P30054 and doc78. The work of the second author was supported in part by the National Science Foundation through Award dms-2111020.

Data availability statement

No new data were created or analysed in this study.

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10.1088/1361-6420/ac6b31